6 مستودعات
Techniques for transforming noise into images using numerical solvers and discretization methods.
Distinct from Monte Carlo Sampling Methods: Existing candidates focus on statistical or MCMC sampling for data analysis, not the iterative denoising process of diffusion models.
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This is a framework for training and sampling diffusion models to generate high-fidelity images, video, and 4D assets. It provides a modular environment for managing generative AI training pipelines, including the handling of datasets, noise sampling, and loss weighting to stabilize the creation of synthetic content. The project features a modular model configuration system that uses YAML-based assembly to define network submodules and conditioners. It also includes a dedicated toolset for AI image watermarking, allowing for the embedding and detection of invisible markers to verify the origi
Provides customizable numerical solvers and discretization methods to generate final outputs from the diffusion model.
Implementation of Denoising Diffusion Probabilistic Model in Pytorch
Generates new data by iteratively applying the learned denoising step from random noise.
This project is a diffusion-based AI art generator and animation framework used to create digital images and motion graphics from text prompts. It functions as a system for producing stylized videos and AI art through iterative diffusion sampling and neural network models. The framework distinguishes itself through specialized tools for 3D depth animation, using depth-map transformations to create spatial movement. It also includes neural style transfer capabilities to apply specific artistic looks, such as watercolor or pixel art, and utilizes optical flow frame blending to reduce flickering
Employs iterative diffusion sampling to generate high-quality digital art from text prompts and visual cues.
This is a classifier-guided diffusion framework for high-fidelity image generation. It implements a cascaded diffusion pipeline that chains a base diffusion model with a dedicated upsampler to progressively increase image resolution in stages, and uses classifier-guided diffusion sampling to steer the reverse diffusion process toward higher-quality outputs. The framework provides tools for training diffusion models from scratch using distributed processes with gradient accumulation, as well as training classifier models that provide gradient-based guidance during sampling. It supports both un
Provides classifier-guided diffusion sampling that injects classifier gradients to steer sample quality toward higher fidelity.
DiffSinger هو مركب صوتي للذكاء الاصطناعي ومولد صوت عصبي مصمم لإنتاج غناء وكلام عالي الدقة. يعمل كنظام تحويل النص إلى كلام وأداة تركيب صوت غنائي قائمة على الانتشار (Diffusion) تحول النص وطبقة الصوت إلى صوت مسموع. يستخدم النظام آلية انتشار ضحلة وتحسين ضوضاء تكراري لإنتاج عروض صوتية واقعية. ويدمج إضافات أخذ عينات متخصصة ومحلات عددية لتسريع الاستدلال وتقليل الوقت المطلوب لتوليد أصوات اصطناعية. يغطي المشروع النمذجة الصوتية، وتركيب مخطط ميل الطيفي (Mel-spectrogram)، وإعادة بناء الموكل العصبي (Neural vocoder) لتحويل النص إلى أشكال موجية صوتية في النطاق الزمني. كما يتضمن قدرات لتحسين الصوت الاصطناعي لتحسين الجودة الصوتية للتسجيلات.
Generates high-fidelity audio by iteratively refining noise into mel-spectrograms using a limited number of sampling steps.
GLIDE is a generative model designed for text-to-image synthesis, image editing, and the contextual filling of masked image regions. It uses a guided diffusion process to transform random noise into high-resolution imagery that aligns with descriptive text prompts. The system provides specialized capabilities for modifying existing visuals, including the ability to alter specific image elements and iteratively refine selected regions through text-driven guidance. It also functions as an inpainting tool, filling missing or masked sections of an image with new content that blends naturally with
Incorporates a separate classifier to guide the diffusion sampling process for better text alignment.