This is a PyTorch-based implementation of diffusion models for synthesizing photorealistic images and video. It provides a framework for text-to-image and text-to-video generation, as well as unconditional image synthesis. The system utilizes a cascading diffusion pipeline to produce high-resolution imagery by passing low-resolution outputs through a sequence of super-resolution models. It also includes capabilities for image inpainting, allowing the reconstruction of masked or missing regions of visual media guided by surrounding context and text prompts. The project includes tools for diff
This is a PyTorch implementation of a text-to-image model designed for synthesizing high-fidelity images from natural language descriptions. It utilizes a diffusion image generator to transform latent embeddings into visual data through an iterative denoising process. The system employs a two-stage latent mapping process, using a CLIP-based latent prior to map text embeddings to image embeddings before decoding them into pixels. It features a cascading diffusion decoder that produces high-resolution imagery by passing low-resolution outputs through a sequence of models at increasing scales.
Latent Diffusion is a framework for high-resolution image synthesis that performs the denoising process within a compressed latent space. It uses variational autoencoders to encode images into a lower-dimensional representation, reducing the computational cost of noise prediction compared to operating on raw pixels. The project enables text-to-image generation by integrating natural language descriptions through cross-attention conditioning. It also supports image inpainting and restoration, filling masked or missing image areas with generated content, and example-based synthesis using retrie
GLIDE is a generative model designed for text-to-image synthesis, image editing, and the contextual filling of masked image regions. It uses a guided diffusion process to transform random noise into high-resolution imagery that aligns with descriptive text prompts. The system provides specialized capabilities for modifying existing visuals, including the ability to alter specific image elements and iteratively refine selected regions through text-driven guidance. It also functions as an inpainting tool, filling missing or masked sections of an image with new content that blends naturally with
本项目是一个专注于生成式扩散模型的教育课程和培训材料集合。它提供了一个课程大纲和实践指南,用于训练、微调和部署能够合成图像、音频和视频的模型。
huggingface/diffusion-models-class 的主要功能包括:Diffusion-Based Media Generation, Classifier-Free Guidance, Diffusion Model Training, Latent Space Compression, Image Inpainting, Cross-Attention Conditioning, Synthetic Media Generators, Concept Fine-Tuning。
huggingface/diffusion-models-class 的开源替代品包括: lucidrains/dalle2-pytorch — This is a PyTorch implementation of a text-to-image model designed for synthesizing high-fidelity images from natural… lucidrains/imagen-pytorch — This is a PyTorch-based implementation of diffusion models for synthesizing photorealistic images and video. It… compvis/latent-diffusion — Latent Diffusion is a framework for high-resolution image synthesis that performs the denoising process within a… openai/glide-text2im — GLIDE is a generative model designed for text-to-image synthesis, image editing, and the contextual filling of masked… openai/improved-diffusion — This project is a diffusion model framework for training and sampling from denoising probabilistic models to generate… open-mmlab/mmagic — mmagic is a multimodal training pipeline and framework for generative AI, focusing on visual synthesis and…