This project is a deep learning framework for AI image super-resolution and facial synthesis. It provides a diffusion model image upscaler and a generative facial image synthesizer capable of transforming low-resolution images into high-resolution outputs using pretrained model weights. The system utilizes iterative diffusion refinement and low-resolution guided sampling to restore fine details and sharpness. It supports both unconditional image generation, where images are created from scratch, and guided resolution enhancement for high-fidelity facial reconstruction. The repository include
This project is a diffusion model framework for training and sampling from denoising probabilistic models to generate images from noise. It functions as a generative image model that creates visual content by iteratively refining random noise into coherent images. The system includes a distributed GPU trainer designed to scale complex neural network architectures across multiple graphics processing units. It also provides an image dataset preprocessor to prepare, scale, and standardize raw image collections for training. The framework covers model training and image generation, utilizing noi
This project is a diffusion model training framework and image synthesis pipeline. It provides the tools necessary to train generative models to learn image data distributions through an iterative denoising process. The framework includes a generative model evaluation tool consisting of automated scripts used to measure the quality and accuracy of produced samples. The system covers model training pipelines and performance evaluation for generative diffusion models.
This is a PyTorch implementation of a text-to-image model designed for synthesizing high-fidelity images from natural language descriptions. It utilizes a diffusion image generator to transform latent embeddings into visual data through an iterative denoising process. The system employs a two-stage latent mapping process, using a CLIP-based latent prior to map text embeddings to image embeddings before decoding them into pixels. It features a cascading diffusion decoder that produces high-resolution imagery by passing low-resolution outputs through a sequence of models at increasing scales.
This is a classifier-guided diffusion framework for high-fidelity image generation. It implements a cascaded diffusion pipeline that chains a base diffusion model with a dedicated upsampler to progressively increase image resolution in stages, and uses classifier-guided diffusion sampling to steer the reverse diffusion process toward higher-quality outputs.
Principalele funcționalități ale openai/guided-diffusion sunt: Classifier-Guided Variants, Image Diffusion Models, Cascaded Pipelines, Classifier-Guided Methods, Diffusion Model Training, Resolution Upscaling, Classifier-Guided Generation, Timestep-Embedded Variants.
Alternativele open-source pentru openai/guided-diffusion includ: janspiry/image-super-resolution-via-iterative-refinement — This project is a deep learning framework for AI image super-resolution and facial synthesis. It provides a diffusion… openai/improved-diffusion — This project is a diffusion model framework for training and sampling from denoising probabilistic models to generate… hojonathanho/diffusion — This project is a diffusion model training framework and image synthesis pipeline. It provides the tools necessary to… lucidrains/dalle2-pytorch — This is a PyTorch implementation of a text-to-image model designed for synthesizing high-fidelity images from natural… lucidrains/denoising-diffusion-pytorch — Implementation of Denoising Diffusion Probabilistic Model in Pytorch. stability-ai/stablecascade — StableCascade is a generative AI system and latent diffusion framework designed for text-to-image synthesis and…