30 open-source projects similar to janspiry/image-super-resolution-via-iterative-refinement, ranked by shared indexed features. Tags may describe platforms or build tools rather than the same primary purpose. Check each project’s use case, license, and deployment requirements before treating it as a replacement.
This project is an AI upscaling framework and deep learning image restorer designed to estimate original source pixels from low-resolution inputs. It functions as a super-resolution reconstruction system that transforms pixelated images into high-resolution versions by restoring high-frequency details and sharpening edges. The system utilizes a convolutional neural network pipeline to analyze pixel data and perform digital image restoration. It employs pixel-shuffle upsampling to rearrange channel dimensions into spatial dimensions, which increases resolution while reducing checkerboard artif
This is a PyTorch implementation of a text-to-image model designed for synthesizing high-fidelity images from natural language descriptions. It utilizes a diffusion image generator to transform latent embeddings into visual data through an iterative denoising process. The system employs a two-stage latent mapping process, using a CLIP-based latent prior to map text embeddings to image embeddings before decoding them into pixels. It features a cascading diffusion decoder that produces high-resolution imagery by passing low-resolution outputs through a sequence of models at increasing scales.
This is a classifier-guided diffusion framework for high-fidelity image generation. It implements a cascaded diffusion pipeline that chains a base diffusion model with a dedicated upsampler to progressively increase image resolution in stages, and uses classifier-guided diffusion sampling to steer the reverse diffusion process toward higher-quality outputs. The framework provides tools for training diffusion models from scratch using distributed processes with gradient accumulation, as well as training classifier models that provide gradient-based guidance during sampling. It supports both un
IOPaint is an AI image editor and Stable Diffusion inpainting tool providing a web interface for removing objects and replacing image content. It utilizes latent diffusion image processing to synthesize high-resolution replacements for erased sections of an image. The project features a specialized AI background remover for isolating subjects and an AI image upscaler that employs super-resolution models for general photos and anime artwork. The software covers a broad range of capabilities including image segmentation for object isolation, face restoration for improving facial details, and t
DiffBIR is a diffusion-based image restoration framework designed for blind image reconstruction. It utilizes generative diffusion priors to recover high-quality images from sources with unknown or complex degradations without requiring explicit degradation models. The system includes specialized models for face restoration, enabling the recovery of facial landmarks, textures, and backgrounds in degraded portraits. To support high-resolution outputs on hardware with limited memory, it employs a tiled image upscaler that divides images into smaller patches during sampling. The framework cover
ESRGAN is a deep learning image restoration framework designed for image super-resolution. It uses a generative adversarial network system to upscale low-resolution images into high-quality versions with sharp visual details and recovered fine textures. The framework implements a perceptual super-resolution model that optimizes the trade-off between perceived visual quality and pixel-level signal-to-noise ratio. It includes weight-interpolation blending to allow for the adjustment of visual sharpness and signal-to-noise ratios by mixing weights from different trained models. The system cover
This project is a deep learning library built for single-image super-resolution and visual enhancement. It provides a framework for training and deploying neural network architectures designed to reconstruct high-resolution images from low-resolution sources, effectively recovering fine details and removing artifacts caused by downscaling or compression. The library distinguishes itself through the implementation of generative adversarial networks and residual block architectures, which work together to improve the realism and clarity of upscaled outputs. It supports training through both pix
This project is a diffusion model framework for training and sampling from denoising probabilistic models to generate images from noise. It functions as a generative image model that creates visual content by iteratively refining random noise into coherent images. The system includes a distributed GPU trainer designed to scale complex neural network architectures across multiple graphics processing units. It also provides an image dataset preprocessor to prepare, scale, and standardize raw image collections for training. The framework covers model training and image generation, utilizing noi
Latent Diffusion is a framework for high-resolution image synthesis that performs the denoising process within a compressed latent space. It uses variational autoencoders to encode images into a lower-dimensional representation, reducing the computational cost of noise prediction compared to operating on raw pixels. The project enables text-to-image generation by integrating natural language descriptions through cross-attention conditioning. It also supports image inpainting and restoration, filling masked or missing image areas with generated content, and example-based synthesis using retrie
This is a PyTorch-based implementation of diffusion models for synthesizing photorealistic images and video. It provides a framework for text-to-image and text-to-video generation, as well as unconditional image synthesis. The system utilizes a cascading diffusion pipeline to produce high-resolution imagery by passing low-resolution outputs through a sequence of super-resolution models. It also includes capabilities for image inpainting, allowing the reconstruction of masked or missing regions of visual media guided by surrounding context and text prompts. The project includes tools for diff
Final2x is an AI image super-resolution tool and neural network inference engine designed to increase image resolution and reconstruct missing details while reducing noise. It functions as a cross-platform image upscaler that executes consistent super-resolution logic across different operating systems. The project serves as a custom model inference engine and upscaling interface, allowing for the import and application of user-defined super-resolution weights and architectures to tailor the visual output of enlarged images. The system utilizes hardware-accelerated processing to offload comp
Pulse is a face image super-resolution tool and self-supervised image enhancer. It functions as a generative model image upsampler and latent space optimization tool designed to increase photo resolution and recover image details. The system differentiates itself by using latent space exploration and spherical constraints to find high-fidelity matches within a generative model. It employs geodesic distance measurement and spherical latent space optimization to regularize representations and maintain parameter radii during the recovery process. The project covers facial image restoration thro
waifu2x-ncnn-vulkan is an AI super-resolution tool and image processor that uses deep learning to increase image resolution and remove visual noise. It is an NCNN-based implementation designed for efficient neural network inference on local hardware. The project utilizes the Vulkan API to provide GPU-accelerated image scaling and noise reduction across diverse graphics hardware. It employs tiled image processing to prevent GPU memory overflow and multi-threaded model loading to reduce initial startup latency. The software covers functional domains including AI image upscaling for maintaining
Clarity-upscaler is an AI image upscaler and enhancement tool that uses deep learning models to increase image resolution and restore visual detail. It functions as a super-resolution inference engine that employs neural networks to predict missing pixels and synthesize high-frequency details from low-resolution sources. The project is delivered as a programmable API, allowing the integration of automated high-resolution image processing and sharpening into external applications and workflows. This interface enables the programmatic upscaling of images to create high-resolution assets. The s
This project is a PyTorch implementation of StyleGAN2, providing a library and research framework for training style-based generative adversarial networks. It serves as a toolkit for high-resolution image synthesis, utilizing competitive minimax optimization to create realistic synthetic visual content. The framework incorporates specialized architectural components such as style-based latent mapping, multi-scale feature modulation, and self-attention layers to improve structural coherence. It distinguishes itself with advanced training stability techniques, including exponential moving avera
BasicSR is a PyTorch-based image restoration toolbox and framework designed for training and deploying deep learning models to upscale, denoise, and deblur images and videos. It serves as a comprehensive system for image super-resolution and video quality restoration, providing the necessary infrastructure to recover fine visual details and increase pixel density. The project distinguishes itself through specialized toolkits for facial image enhancement and high-fidelity face synthesis, as well as a dedicated video quality restoration suite that utilizes deformable convolutions and generative
Sygil-webui is a web interface for Stable Diffusion latent diffusion models, providing a creative suite for text-to-image and text-to-video synthesis. It functions as an image generation tool and a latent diffusion image editor, allowing users to create visuals and video sequences from textual descriptions. The project includes a dedicated model training interface for creating custom textual inversion embeddings, which introduces specific new concepts or styles into the diffusion models. It also features specialized tools for generative image editing, including mask-based inpainting, image-to
Hallo2 is an AI video generation tool and audio-driven portrait animation framework designed to transform static images into speaking videos. It functions as a portrait image animator that synchronizes a single photo with an audio track to produce high-resolution talking head videos. The system includes a distributed animation trainer for fine-tuning deep learning models using custom datasets and distributed computing resources. It employs hierarchical video generation and temporal consistency modeling to produce long-form character animations that remain stable over extended durations. The
This repository provides a complete framework for training generative adversarial networks (GANs) that produce high-resolution photorealistic images, up to 1024 by 1024 pixels. The core technique is progressive layer growth, where both the generator and discriminator networks start training at low resolution and gradually add new layers to model finer details, enabling stable synthesis of large images. The framework includes a high-resolution image generator, an image quality metric evaluator, a latent space interpolation tool for creating smooth transition videos, and a multi-resolution datas
This project is an AI headshot generator and SaaS boilerplate designed to train custom models on uploaded photos to produce professional, high-resolution portraits. It functions as an image generation pipeline and model training orchestrator that manages the end-to-end workflow of processing user images to create stylized avatars. The system includes a credit-based monetization framework that handles payments via automated webhooks. It provides a complete infrastructure for AI-driven services, incorporating user account management and automated email notifications to alert users when their ge
This project is an educational course and collection of training materials focused on generative diffusion models. It provides a curriculum and practical guides for training, fine-tuning, and deploying models capable of synthesizing images, audio, and video. The material covers specific implementation strategies including noise-based synthesis, iterative refinement, and latent space compression. It provides instruction on guiding generative outputs through conditional synthesis and prompt adherence optimization, as well as techniques for image inpainting and text-based editing. The project i
This is a collection of Jupyter notebooks that serve as educational guides for training, fine-tuning, and deploying machine learning models within the Hugging Face ecosystem. The notebooks cover the full lifecycle of model development, from loading and configuring pre-trained transformers to packaging trained models for real-time inference via scalable endpoints. The notebooks demonstrate a range of capabilities including diffusion model training and fine-tuning for image generation and editing, transformer model adaptation for natural language processing tasks, and parameter-efficient fine-t
MochiDiffusion is a local client for Stable Diffusion that functions as an AI image generation studio. It provides a workspace for performing text-to-image, image-to-image, and inpainting tasks, enabling the production of high-resolution images offline using local hardware and neural engine acceleration. The project includes a local model manager for importing, organizing, and converting machine learning models into compatible formats for offline execution. It features a ControlNet integration tool to guide structural composition and spatial layout, alongside a dedicated image upscaler that u
Lama Cleaner is an AI-powered image editing application focused on inpainting, object removal, and generative filling. It provides a suite of tools for erasing unwanted elements from photos and filling the resulting gaps using generative artificial intelligence. The project includes specialized capabilities for image outpainting to extend borders, background removal through object segmentation, and face restoration to fix visual defects. It also features an image upscaler to increase resolution and clarity via super-resolution AI, as well as a Stable Diffusion-based editor for replacing speci
FastVideo is a comprehensive system for accelerated video generation, serving as a video generation inference engine, a video diffusion training framework, and a modular pipeline orchestrator. It provides a distributed transformer optimizer and a distillation toolkit designed to reduce denoising steps and model complexity to increase frame rates. The project distinguishes itself through specialized acceleration techniques, including joint distillation and sparse attention training. It implements low-step video generation and weight quantization to FP8 or FP4 precision to increase throughput a
Srez is a deep learning image super-resolution framework designed to upscale low-resolution images into sharp, high-resolution visual features. It functions as a neural network training tool that employs generative adversarial networks to synthesize realistic image details. The project includes a model evolution visualizer that generates animations and image batches to track visual improvements during the training process. It utilizes a combination of adversarial and L1 loss functions to optimize model weights and supports periodic state checkpointing for recovery and deployment. The system
This PyTorch-based image super-resolution tool provides a deep learning pipeline for upscaling low-resolution images. It utilizes generative adversarial networks to increase pixel density and reconstruct high-resolution image details. The system includes a GAN-based image upscaler and a training pipeline that optimizes neural network weights using paired datasets and custom loss functions. To manage hardware resources, a patch-based image processor splits high-resolution files into smaller segments to prevent memory allocation errors and system crashes. Additional capabilities include the ap
This project is a diffusion model training framework and image synthesis pipeline. It provides the tools necessary to train generative models to learn image data distributions through an iterative denoising process. The framework includes a generative model evaluation tool consisting of automated scripts used to measure the quality and accuracy of produced samples. The system covers model training pipelines and performance evaluation for generative diffusion models.
This project is a comprehensive collection of educational examples and reference implementations for building vision and language models using PyTorch. It serves as a deep learning tutorial covering the end-to-end process of developing neural networks, from initial architecture definition to final production deployment. The repository provides detailed guides on implementing a wide range of domain-specific models, including convolutional neural networks for object detection and segmentation, as well as transformer and recurrent architectures for natural language processing. It emphasizes gene
Super-Gradients is a PyTorch computer vision framework and training library designed for the full lifecycle of vision models. It functions as a deep learning model optimizer and a deployment toolkit for training and fine-tuning models across image classification, object detection, semantic segmentation, and pose estimation tasks. The project provides specific tools for model optimization, including teacher-student knowledge distillation and numerical precision compression to reduce memory and computational requirements. It also includes the implementation of the Yolo-NAS architecture for high